waifu-diffusion
stable-diffusion-webui
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waifu-diffusion | stable-diffusion-webui | |
---|---|---|
28 | 75 | |
1,926 | 2,208 | |
- | - | |
0.0 | 9.8 | |
about 1 year ago | over 1 year ago | |
Python | Python | |
GNU Affero General Public License v3.0 | GNU Affero General Public License v3.0 |
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waifu-diffusion
- What exactly constitutes a new, different model?
- AI: "All Your Horny Belong to Us"
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Cool Japan Diffusion 2.1.1 has been released! 🎉
It's slander to say that the WD team is milking donations in the same way UD is. They already have a public working model which they are continuing to train. The WD team has also implemented features like aspect ratio bucketing into their custom trainer. If they were really milking the project for donations they would have just used the base compvis trainer it was forked from.
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From one of the original DreamBooth authors : Stop using SKS as the initializer word
Use one of the original SD repos, or the code for Waifu Diffusion, or the Smirkingface refactor.
- Stable Diffusion links from around October 4, 2022 that I collected for further processing
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These images of Senko were Generated by AI (Part 2 - Halloween themed)
My model is based off waifu-diffusion.
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What is the name of the AI that created anime art that people keep talking about?
waifu diffusion (free) and NovelAI (paid afaik) is what I know
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Any finetuners (actual finetuning, not Dreambooth) here who can help me find a good learning rate and epoch? At batch size 12/16 (I am using a single A100) with a learning rate of 1e-5/1e-4 I am overfitting (turning all characters into my character) while the style still isn't there yet with ep14...
This is the one I use: https://github.com/harubaru/waifu-diffusion
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Is there a way to tell the AI which artstyle to use?
This is called textual inversion which use many images and a small group of promtps to describe a single element, and there are two methods: with frozen model and with unfrozen model. The first one just creates a new “word” to “guide” the model, in other words, a word that means the same as a chunk of a prompt. The resulting data is a tiny file. The second one (better known dreamboth training) adds new data to de model, so it can copy the character’s characteristics and artist style more effecively resulting in a new model file (cpkt) There are not a single tutorial, so you can find many info on r/StableDiffusion. Also, there is a third method which is based on the original training (many images and its own description) Instruction are here: Training guide I don’t know if this guide is finished though
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Tsukasa at the Beach (Created by a Waifu-Diffusion AI)
The model is a tweaked version of waifu-diffusion using textual inversion. Still not perfect but the AI can create pretty plausible images.
stable-diffusion-webui
- [Stablediffusion] Interface utilisateur Web Diffusion stable
- Generating game concept art
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../../workspace/imgs/txt2img
I am using this one : https://github.com/hlky/stable-diffusion-webui
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How to generate a similar images to an input image *without* a prompt?
Not sure about the script but you can try using this web-ui's img2img tab.
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Enhancing local detail and cohesion by mosaicing
https://github.com/hlky/stable-diffusion-webui now redirects to /sd-webui/stable-diffusion-webui, as though they're the "true" sd-webui.
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reintalled new Hlky update & img2img returns errors (not just where you have to click on mask & back on crop)
As an update, in case anyone else has the issue, after getting some help (thanks u/vedroboev) I installed from here not sure what the difference is, but I got it working.
- Is anyone else unable to use the site?
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Fixing SD images with img2img, am I misunderstanding the concept?
I would pick a version on the github from 8/31 in the stable diffusion repo and then follow step 2a in this guide https://rentry.org/GUItard to transfer the files from this https://github.com/hlky/stable-diffusion-webui/tree/96aba4b36d59803f3817ee60e96a097f54962ae4
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Can't seem to get img2img up and running
This is a bug with the newest UI version. See this.
- Stable Diffusion Img2Img Help
What are some alternatives?
stable-diffusion - This version of CompVis/stable-diffusion features an interactive command-line script that combines text2img and img2img functionality in a "dream bot" style interface, a WebGUI, and multiple features and other enhancements. [Moved to: https://github.com/invoke-ai/InvokeAI]
GFPGAN - GFPGAN aims at developing Practical Algorithms for Real-world Face Restoration.
stable-diffusion-webui - Stable Diffusion web UI
onnx - Open standard for machine learning interoperability
stable-diffusion-ui - Easiest 1-click way to install and use Stable Diffusion on your computer. Provides a browser UI for generating images from text prompts and images. Just enter your text prompt, and see the generated image. [Moved to: https://github.com/easydiffusion/easydiffusion]
diffusers-uncensored - Uncensored fork of diffusers
merge-models - Merges two latent diffusion models at a user-defined ratio
txt2imghd - A port of GOBIG for Stable Diffusion
stable-diffusion-webui - Stable Diffusion web UI [Moved to: https://github.com/Sygil-Dev/sygil-webui]
taming-transformers - Taming Transformers for High-Resolution Image Synthesis
stable-diffusion-gui - Windows GUI for Stable Diffusion