RAVE
diffusers
| RAVE | diffusers | |
|---|---|---|
| 1 | 266 | |
| 213 | 22,646 | |
| 5.6% | 2.8% | |
| 7.2 | 9.9 | |
| about 1 month ago | 6 days ago | |
| Python | Python | |
| MIT License | Apache License 2.0 |
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For example, an activity of 9.0 indicates that a project is amongst the top 10% of the most actively developed projects that we are tracking.
RAVE
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RAVE has been released!
New preprint alert! Introducing RAVE - a zero-shot, lightweight, and fast framework for text-guided video editing, supporting videos of any length utilizing text-to-image pretrained diffusion models. Project Webpage: https://rave-video.github.io ArXiv: https://arxiv.org/abs/2312.04524 More Examples: https://rave-video.github.io/supp/supp.html Code: https://github.com/rehg-lab/RAVE Demo: https://github.com/rehg-lab/RAVE/blob/main/demo_notebook.ipynb Abstract:
diffusers
- StableDiffusionSafetyChecker
- 🧨 diffusers 0.24.0 is out with Kandinsky 3.0, IP Adapters, and others
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What am I missing here? wheres the RND coming from?
I'm missing something about the random factor, from the sample code from https://github.com/huggingface/diffusers/blob/main/README.md
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T2IAdapter+ControlNet at the same time
Hey people, I noticed that combining these two methods in a single forward pass increases the controllability of the generation quite a bit. I was kind of puzzled that sometimes ControlNet yielded better results than T2IAdapter for some cases, and sometimes it was the other way around, so I decided to test both at the same time, and results were quite nice. Some visuals and more motivation here: https://github.com/huggingface/diffusers/issues/5847 And it was already merged here: https://github.com/huggingface/diffusers/pull/5869
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Won't you benchmark me?
Open Parti Prompts: The better way to evaluate diffusion models (repo)
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kohya_ss error. How do I solve this?
You have disabled the safety checker for by passing `safety_checker=None`. Ensure that you abide to the conditions of the Stable Diffusion license and do not expose unfiltered results in services or applications open to the public. Both the diffusers team and Hugging Face strongly recommend to keep the safety filter enabled in all public facing circumstances, disabling it only for use-cases that involve analyzing network behavior or auditing its results. For more information, please have a look at https://github.com/huggingface/diffusers/pull/254 .
- Making a ControlNet inpaint for sdxl
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Stable Diffusion Gets a Major Boost with RTX Acceleration
For developers, TensorRT support also exists for the diffusers library via community pipelines. [1] It's limited, but if you're only supporting a subset of features, it can help.
In general, these insane speed boosts comes at the cost of bleeding edge features.
[1] https://github.com/huggingface/diffusers/blob/28e8d1f6ec82a6...
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Mysterious weights when training UNET
I was training sdxl UNET base model, with the diffusers library, which was going great until around step 210k when the weights suddenly turned back to their original values and stayed that way. I also tried with the ema version, which didn't change at all. I also looked at the tensor's weight values directly which confirmed my suspicions.
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I Made Stable Diffusion XL Smarter by Finetuning It on Bad AI-Generated Images
Merging LoRAs is essentially taking a weighted average of the LoRA adapter weights. It's more common in other UIs.
diffusers is working on a PR for it: https://github.com/huggingface/diffusers/pull/4473
What are some alternatives?
stable-diffusion-webui - Stable Diffusion web UI
riffusion - Stable diffusion for real-time music generation
stable-diffusion - A latent text-to-image diffusion model
lora - Using Low-rank adaptation to quickly fine-tune diffusion models.
invisible-watermark - python library for invisible image watermark (blind image watermark)
automatic - SD.Next: Advanced Implementation of Stable Diffusion and other Diffusion-based generative image models
Dreambooth-Stable-Diffusion - Implementation of Dreambooth (https://arxiv.org/abs/2208.12242) by way of Textual Inversion (https://arxiv.org/abs/2208.01618) for Stable Diffusion (https://arxiv.org/abs/2112.10752). Tweaks focused on training faces, objects, and styles.
sd-webui-additional-networks
sd-webui-modelscope-text2video - Auto1111 extension consisting of implementation of text2video diffusion models (like ModelScope or VideoCrafter) using only Auto1111 webui dependencies [Moved to: https://github.com/deforum-art/sd-webui-text2video]
LoRA - Code for loralib, an implementation of "LoRA: Low-Rank Adaptation of Large Language Models"
gradio - Build and share delightful machine learning apps, all in Python. 🌟 Star to support our work!