ControlNet-v1-1-nightly
diffusers
ControlNet-v1-1-nightly | diffusers | |
---|---|---|
31 | 266 | |
4,349 | 22,763 | |
- | 3.3% | |
8.4 | 9.9 | |
6 months ago | 5 days ago | |
Python | Python | |
- | Apache License 2.0 |
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ControlNet-v1-1-nightly
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Making a ControlNet inpaint for sdxl
1- https://github.com/lllyasviel/ControlNet-v1-1-nightly/issues/89
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AI Yearbook Photos Workflow with Stable Diffusion 1.5 Automatic1111
Install ControlNet and download the models you want to use (canny/depth/openpose should be enough for this): https://github.com/lllyasviel/ControlNet-v1-1-nightly
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can you downgrade Controlnet?
you can find the previous version on their git and if it's a previous version of v1.1 then you probably have to search for the right branch on the new git version and download that
- Could you help me with this problem?
- Controlnet v1.1 Lineart
- Request for current ControlNet information
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AI conceptual massing iterations within a context image with input control sketch
Stable Diffusion: https://huggingface.co/runwayml/stable-diffusion-v1-5 with ControlNet extension: https://github.com/lllyasviel/ControlNet-v1-1-nightly running on Automatic1111 web UI: https://github.com/AUTOMATIC1111/stable-diffusion-webui
- Inpaint Anything (uses "Segment Anything") - Cool A1111 Extension not (yet) on the in App list
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Architectural design using Stable Diffusion and ControlNet
Sure thing, after testing midjourney a bit I found out that yhe quality of images produced is best but you have zero control on over what is produced. The big breakthrough here is ControlNet which is a Stable Diffusion extension that makes you control the initial noise based on image inputs (or at least this is what i understand) more on it here: https://github.com/lllyasviel/ControlNet-v1-1-nightly
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Setting Removed from ControlNET - "Skip img2img processing when using img2img initial image" - why?
https://github.com/lllyasviel/ControlNet-v1-1-nightly/issues/61 it seems get removed as duplicated.
diffusers
- StableDiffusionSafetyChecker
- 🧨 diffusers 0.24.0 is out with Kandinsky 3.0, IP Adapters, and others
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What am I missing here? wheres the RND coming from?
I'm missing something about the random factor, from the sample code from https://github.com/huggingface/diffusers/blob/main/README.md
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T2IAdapter+ControlNet at the same time
Hey people, I noticed that combining these two methods in a single forward pass increases the controllability of the generation quite a bit. I was kind of puzzled that sometimes ControlNet yielded better results than T2IAdapter for some cases, and sometimes it was the other way around, so I decided to test both at the same time, and results were quite nice. Some visuals and more motivation here: https://github.com/huggingface/diffusers/issues/5847 And it was already merged here: https://github.com/huggingface/diffusers/pull/5869
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Won't you benchmark me?
Open Parti Prompts: The better way to evaluate diffusion models (repo)
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kohya_ss error. How do I solve this?
You have disabled the safety checker for by passing `safety_checker=None`. Ensure that you abide to the conditions of the Stable Diffusion license and do not expose unfiltered results in services or applications open to the public. Both the diffusers team and Hugging Face strongly recommend to keep the safety filter enabled in all public facing circumstances, disabling it only for use-cases that involve analyzing network behavior or auditing its results. For more information, please have a look at https://github.com/huggingface/diffusers/pull/254 .
- Making a ControlNet inpaint for sdxl
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Stable Diffusion Gets a Major Boost with RTX Acceleration
For developers, TensorRT support also exists for the diffusers library via community pipelines. [1] It's limited, but if you're only supporting a subset of features, it can help.
In general, these insane speed boosts comes at the cost of bleeding edge features.
[1] https://github.com/huggingface/diffusers/blob/28e8d1f6ec82a6...
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Mysterious weights when training UNET
I was training sdxl UNET base model, with the diffusers library, which was going great until around step 210k when the weights suddenly turned back to their original values and stayed that way. I also tried with the ema version, which didn't change at all. I also looked at the tensor's weight values directly which confirmed my suspicions.
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I Made Stable Diffusion XL Smarter by Finetuning It on Bad AI-Generated Images
Merging LoRAs is essentially taking a weighted average of the LoRA adapter weights. It's more common in other UIs.
diffusers is working on a PR for it: https://github.com/huggingface/diffusers/pull/4473
What are some alternatives?
sd-webui-controlnet - WebUI extension for ControlNet
stable-diffusion-webui - Stable Diffusion web UI
ControlNet - Let us control diffusion models!
stable-diffusion - A latent text-to-image diffusion model
sd-webui-reactor - Fast and Simple Face Swap Extension for StableDiffusion WebUI (A1111 SD WebUI, SD WebUI Forge, SD.Next, Cagliostro)
lora - Using Low-rank adaptation to quickly fine-tune diffusion models.
ControlNet-v1-1-nightly-colab - controlnet v1.1 colab
invisible-watermark - python library for invisible image watermark (blind image watermark)
style2paints - sketch + style = paints :art: (TOG2018/SIGGRAPH2018ASIA)
automatic - SD.Next: Advanced Implementation of Stable Diffusion and other Diffusion-based generative image models
sd-webui-inpaint-anything - Inpaint Anything extension performs stable diffusion inpainting on a browser UI using masks from Segment Anything.
Dreambooth-Stable-Diffusion - Implementation of Dreambooth (https://arxiv.org/abs/2208.12242) by way of Textual Inversion (https://arxiv.org/abs/2208.01618) for Stable Diffusion (https://arxiv.org/abs/2112.10752). Tweaks focused on training faces, objects, and styles.