stable-diffusion-webui-composable-lora VS stable-diffusion-webui

They do, it just needs more work. Stable Diffusion + Latent Couple + Composable LoRA let's you have an Elon model and a Mark model in different parts of the image. No Elon/Mark-ness dominating all the subjects.

https://github.com/opparco/stable-diffusion-webui-two-shothttps://github.com/butaixianran/Stable-Diffusion-Webui-Civitai-Helper https://github.com/opparco/stable-diffusion-webui-composable-lora https://github.com/thomasasfk/sd-webui-aspect-ratio-helper

Possibly with something like the Latent Couple extension + Composable LoRA, cf this video.

I think it works with Regional Prompting in combination with and Composable LoRA (both std. Auto1111-extensions). Never tried myself, though.

So i noticed that Latent Couple and Composable Lora don't work on Vlad Automatic (something to do with those extensions beind made for a outdated version of A1111 webUI).

Composable LoRA Enables using AND keyword(composable diffusion) to limit LoRAs to subprompts. Useful when paired with Latent Couple extension.

Its https://github.com/opparco/stable-diffusion-webui-composable-lora

I tried using composable lora but it is hardcoded to look for the "and", plus it has certain issues that haven't been resolved..

Only one good suggestion can be had here and that is the composable lora extension. I think you need to use it together with the latent couple extension.

Oh this is super easy, just grab the Latent Couple extension and the Composable LoRA extension. Using this you can define multiple sections for different characters, then you can have a prompt/lora unique to each section/character.

* LLaVA model: I'll add more documentation. You are right Llava could not generate images. For image generation I don't have immediate plans, but checkout these projects for local image generation.

- https://diffusionbee.com/

- https://github.com/comfyanonymous/ComfyUI

- https://github.com/AUTOMATIC1111/stable-diffusion-webui

I would love to be able to have a native stable diffusion experience, my rx 580 takes 30s to generate a single image. But it does work after following https://github.com/AUTOMATIC1111/stable-diffusion-webui/wiki...

I got this up and running on my windows machine in short order and I don't even know what stable diffusion is.

But again, it would be nice to have first class support to locally participate in the fun.

In Auto1111, that just uses Image.blend. :)

https://github.com/AUTOMATIC1111/stable-diffusion-webui/blob...

I used SDXL in automatic1111 webui for both images. Now that I think about it, the procedure I described was how I made this one, but the one that looks like an illustration was done in two steps. I used the canny ControlNet as I said for the outer part of the logo to preserve the shape of the fonts, but I had to turn it off for the boot to give SDXL leeway to add detail and make it look more like a boot.

That said, please do learn coding and don't get discouraged when somebody says to learn PyTorch or recommends using a Jupiter notebook with no further information on how to translate the skill into images. I would highly recommend some short term goals. Get your feet wet by taking apart the UIs. The comfy API documentation is here and the A1111 API documentation is here. There is a difference in completeness, welcome to programming. Writing nodes or plugins is also a good way to jump into this world. Custom wildcard logic might be very attractive to you if you aren't the type that want to deal with a nested file structure to simulate logic.

Possibly Unprompted? https://github.com/AUTOMATIC1111/stable-diffusion-webui/discussions/8094

version: v1.6.1  •  python: 3.10.13  •  torch: 2.0.1+cu118  •  xformers: 0.0.20  •  gradio: 3.41.2  •  checkpoint: 6e8d4871f8