Dreambooth-SD-optimized
Stable-Diffusion-Regularization-Images
Dreambooth-SD-optimized | Stable-Diffusion-Regularization-Images | |
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26 | 14 | |
341 | 99 | |
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1.8 | 10.0 | |
over 1 year ago | over 1 year ago | |
Jupyter Notebook | ||
MIT License | - |
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Dreambooth-SD-optimized
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Rtx 4070 Ti which dreambooth could fit
Hey guys im quite new to dream booth i tried the one in stabel diffusion but wasnt really satisfied with the output. Im lloking for a external dreambooth that could be started with anaconda but doesnt need 24 gb Vram i only have 12 Gb. I tried gammagec / Dreambooth-SD-optimizedbut he says you need at least 24 GB
- Best Local SD/Dream Both Combination For Those With 24GB Cards
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Update 1.7.0 of my Windows SD GUI is out! Supports VAE selection, prompt wildcards, even easier DreamBooth training, and tons of quality-of-life improvements. Details in comments.
Using this GitHub https://github.com/gammagec/Dreambooth-SD-optimized from this guide https://pastebin.com/xcFpp9Mr
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Questions about training parameters.
I had pretty good results with 20 images of myself, 200 reg images and 6000 step using https://github.com/gammagec/Dreambooth-SD-optimized.
- [Dreambooth] I changed something about the way Dreambooth training works. Tell me what you think, please.
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First full music video with Deforum 0.5 (single render)
I use Automatic1111 for SD and then Dreambooth Optimized https://github.com/gammagec/Dreambooth-SD-optimized to do custom models.
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How to increase the value of the num_workers?
Gammagec Dreambooth-SD-optimized - https://github.com/gammagec/Dreambooth-SD-optimized
- [Guide] DreamBooth Training with ShivamShrirao's Repo on Windows Locally
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Looking to replicate these kind of effects in stable diffusion. Anyone know what prompts/techniques would be involved? I'd guess they used img2img + ebsynth?
I use this dreambooth repo to train the SD model. https://github.com/gammagec/Dreambooth-SD-optimized Here's the video that shows how to install it in very good detail. https://youtu.be/TwhqmkzdH3s He uses it to train in a face but you can use it to train in a style as well. I suggest taking about 15 or 20 detailed frames of the video and training it in as a style for the class name. You'll have to experiment with how many training steps to take. I suggest doing 1,000 steps at a time and testing out the model. Also don't leave the default "sks", the researchers forgot that that's a common acronym if you know what I mean. Do something like my_style1 so the model doesn't get confused with something else.
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Fewer steps produce more clear images
I'm using this guide: https://www.reddit.com/r/StableDiffusion/comments/xpoexy/yet_another_dreambooth_post_how_to_train_an_image/ to train locally with this repo https://github.com/gammagec/Dreambooth-SD-optimized on windows. Needs a 24gb card though.
Stable-Diffusion-Regularization-Images
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Clarification regularization for Stable Diffusion
However, when I look at regularization dataset that people have created, a lot of them are composed by bad quality AI generated pictures, for instance disfigured humans, or images full or artifacts. For instance, this image of train, or this one of a woman.
- Training Picture Source
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💡 How train with locally with 1.5 Runwayml Inpainting Model?
BTW, you can find the regularization images (ready to use class images) here.
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Regularization images
Have you compared results to using regularization images from an existing repo such as https://github.com/JoePenna/Stable-Diffusion-Regularization-Images
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Comic Diffusion V2. This is a culmination of everything worked towards so far. Trained on 6 styles at the same time, mix and match any number of them to create multiple different unique and consistent styles.
For subjects/people, paste this into the github downloader https://github.com/JoePenna/Stable-Diffusion-Regularization-Images/tree/main/person_ddim
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Good Dreambooth Formula
If you are using person, man or woman as class, you don't need to generate the images as there are a some github repos that have a bunch of them already generated for you to use. Nitrosocke also shared some, check my initial post for the link.
- Custom Model Comparison 1.4 vs 1.5 (something broke)
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What should I do when want better results for a person that was already trained in the sd v1.4 version? Train the model, Dreambooth, or textual inversion embeddings?
I did some experiments with dreambooth training. Overall better results were when I have used 1500 "person" class and about 50 training images. It is vital to have different background and different clothes otherwise it will "bake it" into your token (e.g. same sweater will influence all the rendering with color or pattern as it will be "part of your token"). Now I need to test textual inversion and see the difference.
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Any advice on how to use the dream booth colab with automatic?
As for what kind of images to use I've tried actual photos of people and images generated with Stable Diffusion and I've had pretty good results with both. I also tried using exclusively pictures of the person I'm training for everything and even that worked pretty well. All I can really say is that it seems to pay off if you keep an eye on the framing of your images - if the majority of your reference images cut off the upper 10% of the head for example then your model will tend to also produce images that cut off the upper 10% of the head. Oh, and I haven't tried it myself but this Github repository apparently has a ton of images specifically for use in DreamBooth.
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How are you achieving decent results in DreamBooth? My images look terrible!
I've made sure all my images are only me, and clean images. I have tried using the unsplash regularization images from https://github.com/JoePenna/Stable-Diffusion-Regularization-Images. I've tried generating my own images from SD itself. I've tried 1k, 2k, 3k, 4k steps. I've tried more images of myself and fewer. I've tried using "man", "person", "face" as the class. All of it results in absolute garbage. I get outputs that consistently look like I'm 80 years old or a different ethnicity. Or just wrong... so wrong.
What are some alternatives?
stable-diffusion-webui - Stable Diffusion web UI
SD-Regularization-Images-Style-Dreambooth
Dreambooth-Stable-Diffusion - Implementation of Dreambooth (https://arxiv.org/abs/2208.12242) by way of Textual Inversion (https://arxiv.org/abs/2208.01618) for Stable Diffusion (https://arxiv.org/abs/2112.10752). Tweaks focused on training faces, objects, and styles.
stable-diffusion-ui - Easiest 1-click way to install and use Stable Diffusion on your computer. Provides a browser UI for generating images from text prompts and images. Just enter your text prompt, and see the generated image. [Moved to: https://github.com/easydiffusion/easydiffusion]
Dreambooth-Regularization - All the regs
Dreambooth-Stable-Diffusion - Implementation of Dreambooth (https://arxiv.org/abs/2208.12242) with Stable Diffusion
kohya_ss
Dreambooth-Stable-Diffusion - Implementation of Dreambooth (https://arxiv.org/abs/2208.12242) with Stable Diffusion (tweaks focused on training faces)
fast-stable-diffusion - fast-stable-diffusion + DreamBooth
diffusers - 🤗 Diffusers: State-of-the-art diffusion models for image and audio generation in PyTorch