CrossAttentionControl
diffusers
| CrossAttentionControl | diffusers | |
|---|---|---|
| 11 | 266 | |
| 1,237 | 22,763 | |
| - | 2.8% | |
| 10.0 | 9.9 | |
| over 1 year ago | about 16 hours ago | |
| Jupyter Notebook | Python | |
| MIT License | Apache License 2.0 |
Stars - the number of stars that a project has on GitHub. Growth - month over month growth in stars.
Activity is a relative number indicating how actively a project is being developed. Recent commits have higher weight than older ones.
For example, an activity of 9.0 indicates that a project is amongst the top 10% of the most actively developed projects that we are tracking.
CrossAttentionControl
- "How can I do X?" for image generation.
- The Stable Horde now supports img2img as well as multiple models available at the same time. And we just added SD 1.5
- Is there any way to make Automatic1111 change an image into a different pose/style while keeping the subject of the image in tact?
- Cross Attention Control with Stable Diffusion
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First round of results from the new Cross-Attention paper
Stable Diffusion implementation of Cross Attention Github page (Legend!): https://github.com/bloc97/CrossAttentionControl
- Prompt-to-Prompt Image Editing with Cross Attention Control
- Reproducing the method in 'Prompt-to-Prompt Image Editing with Cross Attention Control' with Stable Diffusion
- Prompt-to-Prompt Image Editing with Cross Attention Control in Stable Diffusion
diffusers
- StableDiffusionSafetyChecker
- 🧨 diffusers 0.24.0 is out with Kandinsky 3.0, IP Adapters, and others
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What am I missing here? wheres the RND coming from?
I'm missing something about the random factor, from the sample code from https://github.com/huggingface/diffusers/blob/main/README.md
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T2IAdapter+ControlNet at the same time
Hey people, I noticed that combining these two methods in a single forward pass increases the controllability of the generation quite a bit. I was kind of puzzled that sometimes ControlNet yielded better results than T2IAdapter for some cases, and sometimes it was the other way around, so I decided to test both at the same time, and results were quite nice. Some visuals and more motivation here: https://github.com/huggingface/diffusers/issues/5847 And it was already merged here: https://github.com/huggingface/diffusers/pull/5869
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Won't you benchmark me?
Open Parti Prompts: The better way to evaluate diffusion models (repo)
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kohya_ss error. How do I solve this?
You have disabled the safety checker for by passing `safety_checker=None`. Ensure that you abide to the conditions of the Stable Diffusion license and do not expose unfiltered results in services or applications open to the public. Both the diffusers team and Hugging Face strongly recommend to keep the safety filter enabled in all public facing circumstances, disabling it only for use-cases that involve analyzing network behavior or auditing its results. For more information, please have a look at https://github.com/huggingface/diffusers/pull/254 .
- Making a ControlNet inpaint for sdxl
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Stable Diffusion Gets a Major Boost with RTX Acceleration
For developers, TensorRT support also exists for the diffusers library via community pipelines. [1] It's limited, but if you're only supporting a subset of features, it can help.
In general, these insane speed boosts comes at the cost of bleeding edge features.
[1] https://github.com/huggingface/diffusers/blob/28e8d1f6ec82a6...
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Mysterious weights when training UNET
I was training sdxl UNET base model, with the diffusers library, which was going great until around step 210k when the weights suddenly turned back to their original values and stayed that way. I also tried with the ema version, which didn't change at all. I also looked at the tensor's weight values directly which confirmed my suspicions.
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I Made Stable Diffusion XL Smarter by Finetuning It on Bad AI-Generated Images
Merging LoRAs is essentially taking a weighted average of the LoRA adapter weights. It's more common in other UIs.
diffusers is working on a PR for it: https://github.com/huggingface/diffusers/pull/4473
What are some alternatives?
stable-diffusion-webui - Stable Diffusion web UI
stable-diffusion
stable-diffusion - A latent text-to-image diffusion model
stable-diffusion-webui - Stable Diffusion web UI [Moved to: https://github.com/Sygil-Dev/sygil-webui]
lora - Using Low-rank adaptation to quickly fine-tune diffusion models.
Magic123 - [ICLR24] Official PyTorch Implementation of Magic123: One Image to High-Quality 3D Object Generation Using Both 2D and 3D Diffusion Priors
invisible-watermark - python library for invisible image watermark (blind image watermark)
nataili - Nataili is a Python library that provides tools for building multimodal AI applications. With its modular design, Nataili makes it easy to use only the tools you need to build custom AI solutions.
automatic - SD.Next: Advanced Implementation of Stable Diffusion and other Diffusion-based generative image models
anima - Turn text into video using Stable Diffusion and Google FILM
Dreambooth-Stable-Diffusion - Implementation of Dreambooth (https://arxiv.org/abs/2208.12242) by way of Textual Inversion (https://arxiv.org/abs/2208.01618) for Stable Diffusion (https://arxiv.org/abs/2112.10752). Tweaks focused on training faces, objects, and styles.