Animator VS stable-diffusion-webui

The Animator v6 script, with its built in "prompt marching" feature, can be a really powerful tool to smoothly move an animation from one subject to the next. To my eye it can produce video sequences that feel gentler and less jarring than what I see from Deforum, while still retaining a lot of SD's latent space trippiness.

Video created using the Animator script :https://github.com/Animator-Anon/Animator for Stable Diffusion using the Automatic 1111 interface : https://github.com/AUTOMATIC1111

Anyone know a decent tutorial for the Animator v6 script for Automatic1111? I've gotten some decent result but feel like there's some things I'm completely missing. For instance, I tried setting a positive prompt via the keyframe dialog, with the hope of getting one image to transition into another (frog into deer? whatever) but I mostly get random outputs. I start with an image I like and it morphs into unrecognizable random junk. It looks cool af, but I don't really understand how to control it.

When I checked the first link for animation, I found this link in the list, looks like a way to start with: https://github.com/Animator-Anon/Animator

sure it's Animator-Anon/Animator: Animator script for SD Web UI (github.com)

But he has "animator" in his img2img scripts which I do not. It's right here: https://github.com/Animator-Anon/Animator/issues

I am quite familiar with python and want to add this to huggingface https://github.com/Animator-Anon/Animator

* LLaVA model: I'll add more documentation. You are right Llava could not generate images. For image generation I don't have immediate plans, but checkout these projects for local image generation.

- https://diffusionbee.com/

- https://github.com/comfyanonymous/ComfyUI

- https://github.com/AUTOMATIC1111/stable-diffusion-webui

I would love to be able to have a native stable diffusion experience, my rx 580 takes 30s to generate a single image. But it does work after following https://github.com/AUTOMATIC1111/stable-diffusion-webui/wiki...

I got this up and running on my windows machine in short order and I don't even know what stable diffusion is.

But again, it would be nice to have first class support to locally participate in the fun.

In Auto1111, that just uses Image.blend. :)

https://github.com/AUTOMATIC1111/stable-diffusion-webui/blob...

I used SDXL in automatic1111 webui for both images. Now that I think about it, the procedure I described was how I made this one, but the one that looks like an illustration was done in two steps. I used the canny ControlNet as I said for the outer part of the logo to preserve the shape of the fonts, but I had to turn it off for the boot to give SDXL leeway to add detail and make it look more like a boot.

That said, please do learn coding and don't get discouraged when somebody says to learn PyTorch or recommends using a Jupiter notebook with no further information on how to translate the skill into images. I would highly recommend some short term goals. Get your feet wet by taking apart the UIs. The comfy API documentation is here and the A1111 API documentation is here. There is a difference in completeness, welcome to programming. Writing nodes or plugins is also a good way to jump into this world. Custom wildcard logic might be very attractive to you if you aren't the type that want to deal with a nested file structure to simulate logic.

Possibly Unprompted? https://github.com/AUTOMATIC1111/stable-diffusion-webui/discussions/8094

version: v1.6.1  •  python: 3.10.13  •  torch: 2.0.1+cu118  •  xformers: 0.0.20  •  gradio: 3.41.2  •  checkpoint: 6e8d4871f8