stability-sdk VS stable-diffusion

I use stable diffusion via stability.ai.

SDXL 0.9 is now available on the Clipdrop by Stability AI platform. Stability AI API and DreamStudio customers will be able to access the model this Monday, 26th June as well as other leading image generating tools like NightCafe.

And for your imaging purposes, you might explore things such as the new StableStudio that Stability AI has. Here's an article about StableStudio. Here are imaging things you can do in StableStudio via its API.

I know there are tons are programs out there but I've made some really cool stuff with this free one https://beta.dreamstudio.ai then rotoscoped it onto a green screen then animated it

Currently Midjourney is the best one, but you need to pay. For low to no cost you could use https://app.leonardo.ai/ or https://beta.dreamstudio.ai/

https://beta.dreamstudio.ai/ (belaj bildoj)

> Of course new SD models are also on the horizon…

SDXL is available at https://beta.dreamstudio.ai/ though they say they're going to release more variants.

I think ControlNet is a lot more interesting than just "better tuned models"; it means there's no line between creating something yourself and asking an AI to do it anymore.

(this is actually generated using the new sdxl beta model in https://beta.dreamstudio.ai )

Download this training project: git clone https://github.com/justinpinkney/stable-diffusion.git

looks like it's using the https://github.com/justinpinkney/stable-diffusion/blob/4ac995b6f663b74dfe65400285e193d4167d259c/scripts/gradio_image_mixer.py to do the bulk of the work meaning the core functionality is built into stable diffusion, seems the UI just isn't built to support it. Their ckpt is here too https://huggingface.co/lambdalabs/image-mixer/tree/main.

Any idea how to make Image Mixer work in this build? On RTX 3060 with 12Gb of memory I get the message:

References: ►Read the full article: https://www.louisbouchard.ai/imagic/ ►Kawar, B., Zada, S., Lang, O., Tov, O., Chang, H., Dekel, T., Mosseri, I. and Irani, M., 2022. Imagic: Text-Based Real Image Editing with Diffusion Models. arXiv preprint arXiv:2210.09276. ► Use it with Stable Diffusion: https://github.com/justinpinkney/stable-diffusion/blob/main/notebooks/imagic.ipynb ►My Newsletter (A new AI application explained weekly to your emails!): https://www.louisbouchard.ai/newsletter/

The notebook: https://github.com/justinpinkney/stable-diffusion/blob/main/notebooks/imagic.ipynb

There's a script for the SD --> Diffusers here: https://github.com/justinpinkney/stable-diffusion/blob/main/scripts/convert_sd_to_diffusers.py

You can start with the github which contains the code: https://github.com/justinpinkney/stable-diffusion

Hmmm, I just double checked the hashes of my local file, what's on huggingface, and what you showed above and they all match. I'm not familiar with that repo, so maybe something weird is going on. I tested it using the original txt2img script in the stable diffusion repo:

*PICK\* (Added Sep. 12, 2022) Web app Stable Diffusion Image Variations by lambdalabs. GitHub repo. Generates variations of an input image without use of a text prompt. Censored.